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yangchen123321/lora-trained-xl-dx_dg-1
lora-trained-xl-dx_dg-1 is a text-to-image model from yangchen123321. Use it when you need an image from a text prompt. It is set up for diffusers. The card lists the license as openrail++.
should probably proofread and complete it, then remove this comment. --
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From the Hugging Face model README
These are yangchen123321/lora-trained-xl-dx_dg-1 LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0.
The weights were trained using DreamBooth.
LoRA for the text encoder was enabled: False.
Special VAE used for training: madebyollin/sdxl-vae-fp16-fix.
You should use 8K, high-voltage power lines, 1 broken frayed cable with twisted and blackened strands, sharp focus on the damaged wire, blurred green mountainous background with distant wind turbines, bright daylight, drone patrol perspective.. to trigger the image generation.
Weights for this model are available in Safetensors format.
Download them in the Files & versions tab.
# TODO: add an example code snippet for running this diffusion pipeline
[TODO: provide examples of latent issues and potential remediations]
[TODO: describe the data used to train the model]