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stablediffusionapi/my-stablediffusion-lora-6484
my-stablediffusion-lora-6484 is a text-to-image model from stablediffusionapi. Use it when you need an image from a text prompt. It is set up for diffusers. The card lists the license as openrail++.
These are LoRA adaption weights for stable-diffusion-v1-5/stable-diffusion-v1-5. The weights were trained on photo of ambika0 man using ModelsLab. LoRA for the text encoder was enabled: False.
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.safetensors25.6 MB · 100%
From the Hugging Face model README
These are LoRA adaption weights for stable-diffusion-v1-5/stable-diffusion-v1-5. The weights were trained on photo of ambika0 man using ModelsLab. LoRA for the text encoder was enabled: False.
!pip install -q transformers accelerate peft diffusers
from diffusers import DiffusionPipeline
import torch
pipe_id = "Lykon/DreamShaper"
pipe = DiffusionPipeline.from_pretrained(pipe_id, torch_dtype=torch.float16).to("cuda")
pipe.load_lora_weights("stablediffusionapi/my-stablediffusion-lora-6484", weight_name="pytorch_lora_weights.safetensors", adapter_name="abc")
prompt = "abc of a hacker with a hoodie"
lora_scale = 0.9
image = pipe(
prompt,
num_inference_steps=30,
cross_attention_kwargs={"scale": 0.9},
generator=torch.manual_seed(0)
).images[0]
image