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stablediffusionapi/my-stablediffusion-lora-2812
my-stablediffusion-lora-2812 is a text-to-image model from stablediffusionapi. Use it when you need an image from a text prompt. It is set up for diffusers. The card lists the license as openrail++.
Downloads · 30 days
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3% of all-time downloads
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From the Hugging Face model README
These are stablediffusionapi/my-stablediffusion-lora-2812 LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0. The weights were trained using Modelslab. LoRA for the text encoder was enabled: False. Special VAE used for training: None.
!pip install -q transformers accelerate peft diffusers
from diffusers import DiffusionPipeline
import torch
pipe_id = "stabilityai/stable-diffusion-xl-base-1.0"
pipe = DiffusionPipeline.from_pretrained(pipe_id, torch_dtype=torch.float16).to("cuda")
pipe.load_lora_weights("stablediffusionapi/my-stablediffusion-lora-2812", weight_name="pytorch_lora_weights.safetensors", adapter_name="abc")
prompt = "abc of a hacker with a hoodie"
lora_scale = 0.9
image = pipe(
prompt,
num_inference_steps=30,
cross_attention_kwargs={"scale": lora_scale},
generator=torch.manual_seed(0)
).images[0]
image
You should use material, mahogany, floor, interior, living room, to trigger the image generation.
Weights for this model are available in Safetensors format. Download them in the Files & versions tab.