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prithivMLmods/SD3.5-Large-Anime-LoRA
SD3.5-Large-Anime-LoRA is a text-to-image model from prithivMLmods. Use it when you need an image from a text prompt. It is set up for diffusers. The card lists the license as creativeml-openrail-m.
Downloads · 30 days
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1% of all-time downloads
All-time downloads
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.safetensors270 MB · 100%
From the Hugging Face model README
The model is still in the training phase. This is not the final version and may contain artifacts and perform poorly in some cases.
prithivMLmods/SD3.5-Large-Anime-LoRA
Image Processing Parameters
| Parameter | Value | Parameter | Value |
|---|---|---|---|
| LR Scheduler | constant | Noise Offset | 0.03 |
| Optimizer | AdamW | Multires Noise Discount | 0.1 |
| Network Dim | 64 | Multires Noise Iterations | 10 |
| Network Alpha | 32 | Repeat & Steps | 25 & 2.7K |
| Epoch | 15 | Save Every N Epochs | 1 |
Labeling: florence2-en(natural language & English)
Total Images Used for Training : 23
import torch
from diffusers import StableDiffusion3Pipeline
pipe = StableDiffusion3Pipeline.from_pretrained("stabilityai/stable-diffusion-3.5-large", torch_dtype=torch.bfloat16)
pipe.load_lora_weights("prithivMLmods/SD3.5-Large-Anime-LoRA", weight_name="SD3.5-Large-Anime-LoRA.safetensors")
pipe.fuse_lora(lora_scale=1.0)
pipe.to("cuda")
prompt = "Man in the style of dark beige and brown, uhd image, youthful protagonists, nonrepresentational photography"
negative_prompt = "(lowres, low quality, worst quality)"
image = pipe(prompt=prompt,
negative_prompt=negative_prompt
num_inference_steps=24,
guidance_scale=4.0,
width=960, height=1280,
).images[0]
image.save(f"example.jpg")


[!WARNING] Trigger words: You should use
Anime 35to trigger the image generation.
Weights for this model are available in Safetensors format.
Download them in the Files & versions tab.