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bannon_style1_LoRA
ilyam0211/bannon_style1_LoRA
bannon_style1_LoRA is a text-to-image model from ilyam0211. Use it when you need an image from a text prompt. It is set up for diffusers. The card lists the license as openrail++.
should probably proofread and complete it, then remove this comment. --
Access
Public
Updated Jun 16, 2025
Repo size
30.3 MB
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At a glance
- Task
- Text-to-Image
- Library
- diffusers
- License
- openrail++
- Access
- Public
- Created
- Jun 16, 2025
- Updated
- Jun 16, 2025
- SHA
- 5f4ba86a
Base models
Model card
From the Hugging Face model README
<!-- This model card has been generated automatically according to the information the training script had access to. You
should probably proofread and complete it, then remove this comment. -->
SDXL LoRA DreamBooth - ilyam0211/bannon_style1_LoRA
<Gallery />Model description
These are ilyam0211/bannon_style1_LoRA LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0.
The weights were trained using DreamBooth.
LoRA for the text encoder was enabled: False.
Special VAE used for training: madebyollin/sdxl-vae-fp16-fix.
Trigger words
You should use art in BANNON style to trigger the image generation.
Download model
Weights for this model are available in Safetensors format.
Download them in the Files & versions tab.
Intended uses & limitations
How to use
# TODO: add an example code snippet for running this diffusion pipeline
Limitations and bias
[TODO: provide examples of latent issues and potential remediations]
Training details
[TODO: describe the data used to train the model]
- Task
- Text-to-Image
- Library
- diffusers
- License
- openrail++
- Created
- Jun 16, 2025
- Updated
- Jun 16, 2025