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austin-k-wang/DiffusionLAIR-SD1.5
DiffusionLAIR-SD1.5 is a text-to-image model from austin-k-wang. Use it when you need an image from a text prompt. It is set up for diffusers. The card lists the license as creativeml-openrail-m.
DiffusionLAIR SD1.5 is a preference-optimized fine-tune of Stable Diffusion v1.5. The UNet was fine-tuned using the Diffusion LAIR objective, while the text encoder, VAE, and other pipeline components were kept frozen.
Downloads · 30 days
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52% of all-time downloads
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.safetensors5.1 GB · 100%
From the Hugging Face model README
DiffusionLAIR SD1.5 is a preference-optimized fine-tune of Stable Diffusion v1.5. The UNet was fine-tuned using the Diffusion LAIR objective, while the text encoder, VAE, and other pipeline components were kept frozen.
runwayml/stable-diffusion-v1-5stable-diffusion-v1-5/stable-diffusion-v1-5451f4fe16113bff5a5d2269ed5ad43b0592e9a14StableDiffusionPipelineimport torch
from diffusers import StableDiffusionPipeline
model_id = "austin-k-wang/DiffusionLAIR-SD1.5"
pipe = StableDiffusionPipeline.from_pretrained(
model_id,
torch_dtype=torch.float16,
use_safetensors=True,
).to("cuda")
generator = torch.Generator(device="cuda").manual_seed(42)
image = pipe(
"A queen with red hair and a green and black dress stands veiled in a highly detailed and elegant digital painting.",
generator=generator,
guidance_scale=7.5,
num_inference_steps=30,
).images[0]
image.save("output.png")