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stable-diffusion-v1-5
Taa97/stable-diffusion-v1-5
stable-diffusion-v1-5 is a text-to-image model from Taa97. Use it when you need an image from a text prompt. It is set up for diffusers. The card lists the license as openrail++.
should probably proofread and complete it, then remove this comment. --
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At a glance
- Task
- Text-to-Image
- Library
- diffusers
- License
- openrail++
- Access
- Public
- Created
- Oct 22, 2024
- Updated
- Oct 22, 2024
- SHA
- e1847b16
Base models
Model card
From the Hugging Face model README
<!-- This model card has been generated automatically according to the information the training script had access to. You
should probably proofread and complete it, then remove this comment. -->
SDXL LoRA DreamBooth - Taa97/stable-diffusion-v1-5
<Gallery />Model description
These are Taa97/stable-diffusion-v1-5 LoRA adaption weights for runwayml/stable-diffusion-v1-5.
The weights were trained using DreamBooth.
LoRA for the text encoder was enabled: False.
Special VAE used for training: madebyollin/sdxl-vae-fp16-fix.
Trigger words
You should use _ to trigger the image generation.
Download model
Weights for this model are available in Safetensors format.
Download them in the Files & versions tab.
Intended uses & limitations
How to use
# TODO: add an example code snippet for running this diffusion pipeline
Limitations and bias
[TODO: provide examples of latent issues and potential remediations]
Training details
[TODO: describe the data used to train the model]
- Task
- Text-to-Image
- Library
- diffusers
- License
- openrail++
- Created
- Oct 22, 2024
- Updated
- Oct 22, 2024