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NathMath/sdxl-0.9-usage
sdxl-0.9-usage is a text-to-image model from NathMath. Use it when you need an image from a text prompt. It is set up for diffusers. The card lists the license as apache-2.0.
- This repo is a tutorial intended to help beginners use the new released model, stable-diffusion-xl-0.9 in ComfyUI, with both the base and refiner models together to achieve a magnificent quality of image generation.
Downloads · 30 days
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Updated Jul 27, 2023
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From the Hugging Face model README
This repo is a tutorial intended to help beginners use the new released model, stable-diffusion-xl-0.9 in ComfyUI, with both the base and refiner models together to achieve a magnificent quality of image generation.
With usable demo interfaces for ComfyUI to use the models (see below)!
After test, it is also useful on SDXL-1.0!
Here is a full tutorial to use stable-diffusion-xl-0.9 FROM ZERO!




When you have loaded into the ComfyUI, it will display this default interface shown below:

Then, find out the "load" button on a floating panel, click, and choose the ".json" file download from THIS REPO, it will quickly load a friendly interface to use the stable-diffusion-xl-0.9 series.

Things will become easy soon.
As labels of all components are shown explicitly in interface, you can find enterpoints and panels by their names.
Then, load "base" and "refiner" models, set relevant parameters, and type your favorite prompts to start your exciting journey.
Note that it will automatically save TWO images of EACH queue, respectively before applying with the refiner, and after.
Positive prompts and Negative prompts in both input box are ought to be identical, but you can freely try whatever you want.
The output quality of the stable-diffusion-xl-0.9 models highly depends on your prompts and parameter settings, so try to adjust them to achieve a better quality. Note that these pretrained models (have not been fine tuned yet) need to be fed with enough details about your PAINTING STYLE, so try to provide sufficient information to ensure the output quality.
Although ComfyUI's model moving technique (move models to VRAM when using and move back to RAM after) can significantly reduce the required VRAM use, please note, possibly, if your models are combined with FP16 and FP32, there could be a round-off error when converting data between these two formats. The main reason why it happens may mainly due to the internal codes, and because of this, you'd better RESTART the interface to reload models from your disk every 200 times, for example, to maintain precise weights and eusure best image quality.
Different from previous diffusion models, this stable-diffusion-xl-0.9 is naturally designed for high resolution image generation, so in order to achieve a decent quality, please use a relatively higher resolution, for recommendation, 1024 * 1024.
This is not the final version of this repo. I will frequently update demos for further applications and advanced functions, so if you are interested, please like it.
Need to be updated...