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claude_style_LoRA
Elizaveta12/claude_style_LoRA
claude_style_LoRA is a text-to-image model from Elizaveta12. Use it when you need an image from a text prompt. It is set up for diffusers. The card lists the license as openrail++.
should probably proofread and complete it, then remove this comment. --
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At a glance
- Task
- Text-to-Image
- Library
- diffusers
- License
- openrail++
- Access
- Public
- Created
- Mar 19, 2026
- Updated
- Mar 20, 2026
- SHA
- 277c32b8
Base models
Model card
From the Hugging Face model README
<!-- This model card has been generated automatically according to the information the training script had access to. You
should probably proofread and complete it, then remove this comment. -->
SDXL LoRA DreamBooth - Elizaveta12/claude_style_LoRA
<Gallery />Model description
These are Elizaveta12/claude_style_LoRA LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0.
The weights were trained using DreamBooth.
LoRA for the text encoder was enabled: False.
Special VAE used for training: madebyollin/sdxl-vae-fp16-fix.
Trigger words
You should use photo collage in Claude Monet style to trigger the image generation.
Download model
Weights for this model are available in Safetensors format.
Download them in the Files & versions tab.
Intended uses & limitations
How to use
# TODO: add an example code snippet for running this diffusion pipeline
Limitations and bias
[TODO: provide examples of latent issues and potential remediations]
Training details
[TODO: describe the data used to train the model]
- Task
- Text-to-Image
- Library
- diffusers
- License
- openrail++
- Fal Ai
- live
- Created
- Mar 19, 2026
- Updated
- Mar 20, 2026