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DKTech/dreambooth-test-1
dreambooth-test-1 is a text-to-image model from DKTech. Use it when you need an image from a text prompt. It is set up for diffusers. The card lists the license as openrail++.
should probably proofread and complete it, then remove this comment. --
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From the Hugging Face model README
These are DKTech/dreambooth-test-1 LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0.
The weights were trained using DreamBooth.
LoRA for the text encoder was enabled: False.
Special VAE used for training: madebyollin/sdxl-vae-fp16-fix.
You should use a photo of sks dog to trigger the image generation.
Weights for this model are available in Safetensors format.
Download them in the Files & versions tab.
Set up the environment on command-line / terminal.
# Create and activate conda environment
conda create –name dreambooth python=3.10
conda activate dreambooth
# Install ipykernel (needed only if you want to run the inference inside a jupyter-notebook)
conda install -c anaconda ipykernel
python -m ipykernel install --user --name=dreambooth
# Clone and install diffusers package
git clone https://github.com/huggingface/diffusers
cd diffusers
pip install -e .
# Browse to examples/dreambooth directory in the diffusers installation directory
cd examples/dreambooth
# Install dreambooth sdxl training dependencies
pip install -r requirements_sdxl.txt
Run the inference in Python.
from huggingface_hub.repocard import RepoCard
from diffusers import DiffusionPipeline
import torch
lora_model_id = "DKTech/dreambooth-test-1"
card = RepoCard.load(lora_model_id)
base_model_id = card.data.to_dict()["base_model"]
pipe = DiffusionPipeline.from_pretrained(base_model_id, torch_dtype=torch.float16)
pipe = pipe.to("cuda")
pipe.load_lora_weights(lora_model_id)
image = pipe("A picture of an elephant that looks like a dog.", num_inference_steps=25).images[0]
image.save("my_image.png")
This model was created by fine tuning the original stable diffusion model based on the instructions here- https://github.com/huggingface/diffusers/blob/main/examples/dreambooth/README_sdxl.md Various other base models (other than stable diffusion) can also be fine tuned using DreamBooth. For example, some discussion on fine tuning Playground 2.5 model can be found here- https://github.com/huggingface/diffusers/pull/7126
[TODO: provide examples of latent issues and potential remediations]
[TODO: describe the data used to train the model]